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style_lr_10 – AI Model by kyoungmiin | AlphaNeural AI
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style_lr_10
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diffusers
text-to-image
diffusers-training
lora
sd-lora
stable-diffusion-xl
stable-diffusion-xl-diffusers
stabilityai/stable-diffusion-xl-base-1.0
adapter
openrail++
us
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SDXL LoRA DreamBooth - kyoungmiin/style_lr_10
Model description
These are kyoungmiin/style_lr_10 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using
DreamBooth
.
LoRA for the text encoder was enabled: False.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
Trigger words
You should use sks to trigger the image generation.
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Weights for this model are available in Safetensors format.
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them in the Files & versions tab.
Intended uses & limitations
How to use
# TODO: add an example code snippet for running this diffusion pipeline
Limitations and bias
[TODO: provide examples of latent issues and potential remediations]
Training details
[TODO: describe the data used to train the model]
